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dribnet /dollarpeople:abe86f1f
Input schema
The fields you can use to run this model with an API. If you don’t give a value for a field its default value will be used.
| Field | Type | Default value | Description |
|---|---|---|---|
| prompt |
string
|
Prompt for generated image. If you include the `trigger_word` used in the training process you are more likely to activate the trained object, style, or concept in the resulting image.
|
|
| image |
string
|
Input image for img2img or inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.
|
|
| mask |
string
|
Input mask for inpainting mode. Black areas will be preserved, white areas will be inpainted. Must be provided along with 'image' for inpainting mode.
|
|
| aspect_ratio |
string
(enum)
|
1:1
Options: 1:1, 16:9, 21:9, 3:2, 2:3, 4:5, 5:4, 3:4, 4:3, 9:16, 9:21, custom |
Aspect ratio for the generated image in text-to-image mode. The size will always be 1 megapixel, i.e. 1024x1024 if aspect ratio is 1:1. To use arbitrary width and height, set aspect ratio to 'custom'. Note: Ignored in img2img and inpainting modes.
|
| width |
integer
|
Min: 256 Max: 1440 |
Width of the generated image in text-to-image mode. Only used when aspect_ratio=custom. Must be a multiple of 16 (if it's not, it will be rounded to nearest multiple of 16). Note: Ignored in img2img and inpainting modes.
|
| height |
integer
|
Min: 256 Max: 1440 |
Height of the generated image in text-to-image mode. Only used when aspect_ratio=custom. Must be a multiple of 16 (if it's not, it will be rounded to nearest multiple of 16). Note: Ignored in img2img and inpainting modes.
|
| num_outputs |
integer
|
1
Min: 1 Max: 4 |
Number of images to output.
|
| lora_scale |
number
|
1
Min: -1 Max: 2 |
Determines how strongly the main LoRA should be applied. Sane results between 0 and 1.
|
| num_inference_steps |
integer
|
28
Min: 1 Max: 50 |
Number of inference steps. More steps can give more detailed images, but take longer.
|
| model |
string
(enum)
|
dev
Options: dev, schnell |
Which model to run inferences with. The dev model needs around 28 steps but the schnell model only needs around 4 steps.
|
| guidance_scale |
number
|
3.5
Max: 10 |
Guidance scale for the diffusion process. Lower values can give more realistic images. Good values to try are 2, 2.5, 3 and 3.5
|
| prompt_strength |
number
|
0.8
Max: 1 |
Prompt strength when using img2img / inpaint. 1.0 corresponds to full destruction of information in image
|
| seed |
integer
|
Random seed. Set for reproducible generation.
|
|
| extra_lora |
string
|
Combine this fine-tune with another LoRA. Supports Replicate models in the format <owner>/<username> or <owner>/<username>/<version>, HuggingFace URLs in the format huggingface.co/<owner>/<model-name>, CivitAI URLs in the format civitai.com/models/<id>[/<model-name>], or arbitrary .safetensors URLs from the Internet. For example, 'fofr/flux-pixar-cars'
|
|
| extra_lora_scale |
number
|
1
Min: -1 Max: 2 |
Determines how strongly the extra LoRA should be applied.
|
| output_format |
string
(enum)
|
webp
Options: webp, jpg, png |
Format of the output images.
|
| output_quality |
integer
|
90
Max: 100 |
Quality when saving the output images, from 0 to 100. 100 is best quality, 0 is lowest quality. Not relevant for .png outputs
|
| disable_safety_checker |
boolean
|
False
|
Disable safety checker for generated images.
|
Output schema
The shape of the response you’ll get when you run this model with an API.
{'items': {'format': 'uri', 'type': 'string'},
'title': 'Output',
'type': 'array'}