FLUX.1-Dev LoRA Training (with 2x GPUs) by Huggingface Diffusers
Shiba stable diffusion model
bigcode/tiny_starcoder_py
lmsys/vicuna-13b-v1.3
lmsys/vicuna-7b-v1.3
Salesforce/codegen2-1B
Salesforce/xgen-7b-8k-base
Realistic Vision V3.0 with VAE
Realistic Vision V4.0
CLIP Interrogator (for faster inference)
A working wsrglow model
RiversHaveWings Stable Diffusion Upscaler
Real-ESRGAN with optional face correction and adjustable upscale (for larger images)
Animate Your Personalized Text-to-Image Diffusion Models
Segments an audio recording based on who is speaking (on A100)
Meta's Llama 2 7b Chat - GPTQ
Meta's Llama 2 13b Chat - GPTQ
Stability AI's FreeWilly2
Implementation of Realistic Vision v5.1 with VAE
Practical face restoration algorithm for *old photos* or *AI-generated faces* (for larger images)
SDXL v1.0 - A text-to-image generative AI model that creates beautiful images
This model is not yet booted but ready for API calls. Your first API call will boot the model and may take longer, but after that subsequent responses will be fast.
This model runs on 2x A100 (80GB).
