Collections

Generate images

These models generate images from text prompts. Many of these models are based on Stable Diffusion and FLUX.1.

Read our guide to learn more about using Stable Diffusion.

  • Text-to-image - Convert text prompts to photorealistic images. Useful for quickly visualizing concepts
  • Control over style - Adjust image properties like lighting and texture via prompts
  • In-painting - Expand, edit, or refine images by filling in missing regions

Our Picks

Best overall image generation model: black-forest-labs/flux-pro

The best overall image generation model is black-forest-labs/flux-pro. It offers state-of-the-art performance in prompt following, visual quality, image detail, and output diversity. For more information about how to use FLUX.1, read our blog about FLUX.1.

Best ComfyUI model: fofr/any-comfyui-workflow

If you’re a fan of ComfyUI, you can export any of your favorite ComfyUI workflows to JSON and run them on Replicate using the fofr/any-comfyui-workflow model. For more information, check out our detailed guide to using ComfyUI.

Best fast image generation model: black-forest-labs/flux-schnell

The best-looking fast image generation model is black-forest-labs/flux-schnell, which can generate high-quality images in roughly 1 second. The fastest image generation model is fofr/latent-consistency-model which will generate an image in 0.6 seconds.

Best fine-tunes

Make sure to check out our SDXL fine-tunes collection, which includes all publicly available SDXL fine-tunes hosted on Replicate. This collection should help you get a feel for the sorts of things you can do with fine-tuning.

Recommended models

bytedance / sdxl-lightning-4step

SDXL-Lightning by ByteDance: a fast text-to-image model that makes high-quality images in 4 steps

394.1M runs

stability-ai / stable-diffusion

A latent text-to-image diffusion model capable of generating photo-realistic images given any text input

108.8M runs

stability-ai / sdxl

A text-to-image generative AI model that creates beautiful images

64.1M runs

black-forest-labs / flux-schnell

The fastest image generation model tailored for local development and personal use

30.5M runs

stability-ai / stable-diffusion-inpainting

Fill in masked parts of images with Stable Diffusion

18.1M runs

ai-forever / kandinsky-2.2

multilingual text2image latent diffusion model

10M runs

datacte / proteus-v0.2

Proteus v0.2 shows subtle yet significant improvements over Version 0.1. It demonstrates enhanced prompt understanding that surpasses MJ6, while also approaching its stylistic capabilities.

6.9M runs

ai-forever / kandinsky-2

text2img model trained on LAION HighRes and fine-tuned on internal datasets

6.2M runs

fofr / sdxl-emoji

An SDXL fine-tune based on Apple Emojis

5.6M runs

tstramer / material-diffusion

Stable diffusion fork for generating tileable outputs using v1.5 model

2.1M runs

playgroundai / playground-v2.5-1024px-aesthetic

Playground v2.5 is the state-of-the-art open-source model in aesthetic quality

1.3M runs

black-forest-labs / flux-dev

A 12 billion parameter rectified flow transformer capable of generating images from text descriptions

1.3M runs

datacte / proteus-v0.3

ProteusV0.3: The Anime Update

1.3M runs

fofr / latent-consistency-model

Super-fast, 0.6s per image. LCM with img2img, large batching and canny controlnet

1.1M runs

lucataco / ssd-1b

Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities

955.3K runs

batouresearch / sdxl-controlnet-lora

'''Last update: Now supports img2img.''' SDXL Canny controlnet with LoRA support.

621.8K runs

fofr / realvisxl-v3-multi-controlnet-lora

RealVisXl V3 with multi-controlnet, lora loading, img2img, inpainting

618.7K runs

fofr / any-comfyui-workflow

Run any ComfyUI workflow. Guide: https://github.com/fofr/cog-comfyui

580.3K runs

fofr / sticker-maker

Make stickers with AI. Generates graphics with transparent backgrounds.

491.6K runs

lucataco / realvisxl2-lcm

RealvisXL-v2.0 with LCM LoRA - requires fewer steps (4 to 8 instead of the original 40 to 50)

288.8K runs

lucataco / realvisxl-v2.0

Implementation of SDXL RealVisXL_V2.0

277.6K runs

fofr / sdxl-multi-controlnet-lora

Multi-controlnet, lora loading, img2img, inpainting

200.2K runs

lucataco / dreamshaper-xl-turbo

DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.

171.7K runs

lucataco / open-dalle-v1.1

A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension

102.8K runs

adirik / realvisxl-v3.0-turbo

Photorealism with RealVisXL V3.0 Turbo based on SDXL

88.8K runs

ai-forever / kandinsky-2-1

Kandinsky 2.1 Diffusion Model

83.7K runs

nightmareai / disco-diffusion

Generate images using a variety of techniques - Powered by Discoart

64.2K runs

lucataco / pixart-xl-2

PixArt-Alpha 1024px is a transformer-based text-to-image diffusion system trained on text embeddings from T5

46.6K runs

adirik / realvisxl-v4.0

Photorealism with RealVisXL V4.0

40.4K runs

lucataco / realistic-vision-v5

Realistic Vision v5.0 with VAE

18.7K runs