Collections

Generate images

These models generate images from text prompts. Many of these models are based on Stable Diffusion and FLUX.1.

Read our guide to learn more about using Stable Diffusion.

Our Picks

Best overall image generation model: black-forest-labs/flux-1.1-pro

The best overall image generation model is black-forest-labs/flux-1.1-pro. It offers state-of-the-art performance in prompt following, visual quality, image detail, and output diversity. For more information about how to use FLUX.1, read our blog about FLUX 1.1 pro and check out our collection of the FLUX family of models.

Best fast image generation model: black-forest-labs/flux-schnell

The smallest of the FLUX family of models, black-forest-labs/flux-schnell can generate high-quality images in roughly 1 second.

Best model for generating images with text: ideogram-ai/ideogram-v2

Ideogram models are strong in many areas, but they’re especially known for their ability to generate realistic, legible text. Ideogram v2 also has powerful inpainting features. For more on inpainting with an API, see our blog on Ideogram v2. Or try the live demo of inpainting with Ideogram right away.

Best model for generating images with SVGs: recraft-ai/recraft-v3-svg

The Recraft V3 SVG model is the first major text-to-image model with the ability to generate high quality SVG images including logotypes, and icons. The model supports a wide list of styles.

Best ComfyUI model: fofr/any-comfyui-workflow

If you’re a fan of ComfyUI, you can export any of your favorite ComfyUI workflows to JSON and run them on Replicate using the fofr/any-comfyui-workflow model. For more information, check out our detailed guide to using ComfyUI.

Best fine-tunes

Make sure to check out our SDXL fine-tunes collection, which includes all publicly available SDXL fine-tunes hosted on Replicate, and our FLUX fine-tunes collection as well. This collection should help you get a feel for the sorts of things you can do with fine-tuning.

Recommended models

bytedance / sdxl-lightning-4step

SDXL-Lightning by ByteDance: a fast text-to-image model that makes high-quality images in 4 steps

568.1M runs

black-forest-labs / flux-schnell

The fastest image generation model tailored for local development and personal use

126M runs

stability-ai / stable-diffusion

A latent text-to-image diffusion model capable of generating photo-realistic images given any text input

109.6M runs

stability-ai / sdxl

A text-to-image generative AI model that creates beautiful images

70.4M runs

jagilley / controlnet-scribble

Generate detailed images from scribbled drawings

38.1M runs

stability-ai / stable-diffusion-inpainting

Fill in masked parts of images with Stable Diffusion

18.6M runs

ai-forever / kandinsky-2.2

multilingual text2image latent diffusion model

10M runs

datacte / proteus-v0.2

Proteus v0.2 shows subtle yet significant improvements over Version 0.1. It demonstrates enhanced prompt understanding that surpasses MJ6, while also approaching its stylistic capabilities.

8.1M runs

fofr / sdxl-emoji

An SDXL fine-tune based on Apple Emojis

7.2M runs

black-forest-labs / flux-pro

State-of-the-art image generation with top of the line prompt following, visual quality, image detail and output diversity.

7M runs

ai-forever / kandinsky-2

text2img model trained on LAION HighRes and fine-tuned on internal datasets

6.2M runs

black-forest-labs / flux-dev

A 12 billion parameter rectified flow transformer capable of generating images from text descriptions

4.8M runs

black-forest-labs / flux-1.1-pro

Faster, better FLUX Pro. Text-to-image model with excellent image quality, prompt adherence, and output diversity.

4.7M runs

datacte / proteus-v0.3

ProteusV0.3: The Anime Update

2.2M runs

tstramer / material-diffusion

Stable diffusion fork for generating tileable outputs using v1.5 model

2.2M runs

playgroundai / playground-v2.5-1024px-aesthetic

Playground v2.5 is the state-of-the-art open-source model in aesthetic quality

1.9M runs

fofr / latent-consistency-model

Super-fast, 0.6s per image. LCM with img2img, large batching and canny controlnet

1.2M runs

lucataco / ssd-1b

Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities

989.3K runs

fofr / realvisxl-v3-multi-controlnet-lora

RealVisXl V3 with multi-controlnet, lora loading, img2img, inpainting

925.4K runs

batouresearch / sdxl-controlnet-lora

'''Last update: Now supports img2img.''' SDXL Canny controlnet with LoRA support.

728.3K runs

fofr / any-comfyui-workflow

Run any ComfyUI workflow. Guide: https://github.com/fofr/cog-comfyui

710.5K runs

black-forest-labs / flux-1.1-pro-ultra

FLUX1.1 [pro] in ultra and raw modes. Images are up to 4 megapixels. Use raw mode for realism.

692.5K runs

fofr / sticker-maker

Make stickers with AI. Generates graphics with transparent backgrounds.

557.2K runs

lucataco / realvisxl2-lcm

RealvisXL-v2.0 with LCM LoRA - requires fewer steps (4 to 8 instead of the original 40 to 50)

290.1K runs

lucataco / realvisxl-v2.0

Implementation of SDXL RealVisXL_V2.0

280.1K runs

fofr / sdxl-multi-controlnet-lora

Multi-controlnet, lora loading, img2img, inpainting

204.4K runs

lucataco / dreamshaper-xl-turbo

DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.

185.2K runs

lucataco / open-dalle-v1.1

A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension

121.5K runs

adirik / realvisxl-v3.0-turbo

Photorealism with RealVisXL V3.0 Turbo based on SDXL

109.6K runs

stability-ai / stable-diffusion-3.5-large

A text-to-image model that generates high-resolution images with fine details. It supports various artistic styles and produces diverse outputs from the same prompt, thanks to Query-Key Normalization.

101.9K runs

ai-forever / kandinsky-2-1

Kandinsky 2.1 Diffusion Model

84.1K runs

nightmareai / disco-diffusion

Generate images using a variety of techniques - Powered by Discoart

64.3K runs

ideogram-ai / ideogram-v2-turbo

A fast image model with state of the art inpainting, prompt comprehension and text rendering.

52.5K runs

lucataco / pixart-xl-2

PixArt-Alpha 1024px is a transformer-based text-to-image diffusion system trained on text embeddings from T5

51.4K runs

adirik / realvisxl-v4.0

Photorealism with RealVisXL V4.0

42.9K runs

ideogram-ai / ideogram-v2

An excellent image model with state of the art inpainting, prompt comprehension and text rendering

36.6K runs

lucataco / realistic-vision-v5

Realistic Vision v5.0 with VAE

26.8K runs

stability-ai / stable-diffusion-3.5-large-turbo

A text-to-image model that generates high-resolution images with fine details. It supports various artistic styles and produces diverse outputs from the same prompt, with a focus on fewer inference steps

11.3K runs

stability-ai / stable-diffusion-3.5-medium

2.5 billion parameter image model with improved MMDiT-X architecture

6K runs