Generate images
These models generate images from text prompts. Many of these models are based on Stable Diffusion and FLUX.1.
Read our guide to learn more about using Stable Diffusion.
Our Picks
Best overall image generation model: black-forest-labs/flux-1.1-pro
The best overall image generation model is black-forest-labs/flux-1.1-pro. It offers state-of-the-art performance in prompt following, visual quality, image detail, and output diversity. For more information about how to use FLUX.1, read our blog about FLUX 1.1 pro and check out our collection of the FLUX family of models.
Best fast image generation model: black-forest-labs/flux-schnell
The smallest of the FLUX family of models, black-forest-labs/flux-schnell can generate high-quality images in roughly 1 second.
Best model for generating images with text: ideogram-ai/ideogram-v2
Ideogram models are strong in many areas, but they’re especially known for their ability to generate realistic, legible text. Ideogram v2 also has powerful inpainting features. For more on inpainting with an API, see our blog on Ideogram v2. Or try the live demo of inpainting with Ideogram right away.
Best model for generating images with SVGs: recraft-ai/recraft-v3-svg
The Recraft V3 SVG model is the first major text-to-image model with the ability to generate high quality SVG images including logotypes, and icons. The model supports a wide list of styles.
Best ComfyUI model: fofr/any-comfyui-workflow
If you’re a fan of ComfyUI, you can export any of your favorite ComfyUI workflows to JSON and run them on Replicate using the fofr/any-comfyui-workflow model. For more information, check out our detailed guide to using ComfyUI.
Best fine-tunes
Make sure to check out our SDXL fine-tunes collection, which includes all publicly available SDXL fine-tunes hosted on Replicate, and our FLUX fine-tunes collection as well. This collection should help you get a feel for the sorts of things you can do with fine-tuning.
Recommended models
bytedance / sdxl-lightning-4step
SDXL-Lightning by ByteDance: a fast text-to-image model that makes high-quality images in 4 steps
black-forest-labs / flux-schnell
The fastest image generation model tailored for local development and personal use
stability-ai / stable-diffusion
A latent text-to-image diffusion model capable of generating photo-realistic images given any text input
stability-ai / sdxl
A text-to-image generative AI model that creates beautiful images
jagilley / controlnet-scribble
Generate detailed images from scribbled drawings
stability-ai / stable-diffusion-inpainting
Fill in masked parts of images with Stable Diffusion
ai-forever / kandinsky-2.2
multilingual text2image latent diffusion model
datacte / proteus-v0.2
Proteus v0.2 shows subtle yet significant improvements over Version 0.1. It demonstrates enhanced prompt understanding that surpasses MJ6, while also approaching its stylistic capabilities.
fofr / sdxl-emoji
An SDXL fine-tune based on Apple Emojis
black-forest-labs / flux-pro
State-of-the-art image generation with top of the line prompt following, visual quality, image detail and output diversity.
ai-forever / kandinsky-2
text2img model trained on LAION HighRes and fine-tuned on internal datasets
black-forest-labs / flux-dev
A 12 billion parameter rectified flow transformer capable of generating images from text descriptions
black-forest-labs / flux-1.1-pro
Faster, better FLUX Pro. Text-to-image model with excellent image quality, prompt adherence, and output diversity.
datacte / proteus-v0.3
ProteusV0.3: The Anime Update
tstramer / material-diffusion
Stable diffusion fork for generating tileable outputs using v1.5 model
playgroundai / playground-v2.5-1024px-aesthetic
Playground v2.5 is the state-of-the-art open-source model in aesthetic quality
fofr / latent-consistency-model
Super-fast, 0.6s per image. LCM with img2img, large batching and canny controlnet
lucataco / ssd-1b
Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities
fofr / realvisxl-v3-multi-controlnet-lora
RealVisXl V3 with multi-controlnet, lora loading, img2img, inpainting
batouresearch / sdxl-controlnet-lora
'''Last update: Now supports img2img.''' SDXL Canny controlnet with LoRA support.
fofr / any-comfyui-workflow
Run any ComfyUI workflow. Guide: https://github.com/fofr/cog-comfyui
black-forest-labs / flux-1.1-pro-ultra
FLUX1.1 [pro] in ultra and raw modes. Images are up to 4 megapixels. Use raw mode for realism.
fofr / sticker-maker
Make stickers with AI. Generates graphics with transparent backgrounds.
lucataco / realvisxl2-lcm
RealvisXL-v2.0 with LCM LoRA - requires fewer steps (4 to 8 instead of the original 40 to 50)
lucataco / realvisxl-v2.0
Implementation of SDXL RealVisXL_V2.0
fofr / sdxl-multi-controlnet-lora
Multi-controlnet, lora loading, img2img, inpainting
lucataco / dreamshaper-xl-turbo
DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.
lucataco / open-dalle-v1.1
A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension
adirik / realvisxl-v3.0-turbo
Photorealism with RealVisXL V3.0 Turbo based on SDXL
stability-ai / stable-diffusion-3.5-large
A text-to-image model that generates high-resolution images with fine details. It supports various artistic styles and produces diverse outputs from the same prompt, thanks to Query-Key Normalization.
ai-forever / kandinsky-2-1
Kandinsky 2.1 Diffusion Model
nightmareai / disco-diffusion
Generate images using a variety of techniques - Powered by Discoart
ideogram-ai / ideogram-v2-turbo
A fast image model with state of the art inpainting, prompt comprehension and text rendering.
lucataco / pixart-xl-2
PixArt-Alpha 1024px is a transformer-based text-to-image diffusion system trained on text embeddings from T5
adirik / realvisxl-v4.0
Photorealism with RealVisXL V4.0
ideogram-ai / ideogram-v2
An excellent image model with state of the art inpainting, prompt comprehension and text rendering
lucataco / realistic-vision-v5
Realistic Vision v5.0 with VAE
stability-ai / stable-diffusion-3.5-large-turbo
A text-to-image model that generates high-resolution images with fine details. It supports various artistic styles and produces diverse outputs from the same prompt, with a focus on fewer inference steps
stability-ai / stable-diffusion-3.5-medium
2.5 billion parameter image model with improved MMDiT-X architecture