Explore
Fine-tune FLUX fast
Customize FLUX.1 [dev] with the fast FLUX trainer on Replicate
Train the model to recognize and generate new concepts using a small set of example images, for specific styles, characters, or objects. It's fast (under 2 minutes), cheap (under $2), and gives you a warm, runnable model plus LoRA weights to download.
Featured models
luma / reframe-video
Change the aspect ratio of any video up to 30 seconds long, outputs will be 720p
google / imagen-4-fast
Use this fast version of Imagen 4 when speed and cost are more important than quality
google / imagen-4-ultra
Use this ultra version of Imagen 4 when quality matters more than speed and cost
google / imagen-4
Google's Imagen 4 flagship model
replicate / fast-flux-trainer
Train subjects or styles faster than ever
google / veo-3
Sound on: Google’s flagship Veo 3 text to video model, with audio
black-forest-labs / flux-kontext-pro
A state-of-the-art text-based image editing model that delivers high-quality outputs with excellent prompt following and consistent results for transforming images through natural language
black-forest-labs / flux-kontext-max
A premium text-based image editing model that delivers maximum performance and improved typography generation for transforming images through natural language prompts
ideogram-ai / ideogram-v3-turbo
Turbo is the fastest and cheapest Ideogram v3. v3 creates images with stunning realism, creative designs, and consistent styles
Official models
Official models are always on, maintained, and have predictable pricing.
I want to…
Generate images
Models that generate images from text prompts
Make videos with Wan2.1
Generate videos with Wan2.1, the fastest and highest quality open-source video generation model.
Generate videos
Models that create and edit videos
Caption images
Models that generate text from images
Transcribe speech
Models that convert speech to text
Use the FLUX family of models
The FLUX family of text-to-image models from Black Forest Labs
Remove backgrounds
Models that remove backgrounds from images and videos
Restore images
Models that improve or restore images by deblurring, colorization, and removing noise
Caption videos
Models that generate text from videos
Edit images
Tools for manipulating images.
Use a face to make images
Make realistic images of people instantly
Get embeddings
Models that generate embeddings from inputs
Generate speech
Convert text to speech
Generate music
Models to generate and modify music
Generate text
Models that can understand and generate text
Use handy tools
Toolbelt-type models for videos and images.
Upscale images
Upscaling models that create high-quality images from low-quality images
Use official models
Official models are always on, maintained, and have predictable pricing.
Enhance videos
Models that enhance videos with super-resolution, sound effects, motion capture and other useful production effects.
Detect objects
Models that detect or segment objects in images and videos.
Sing with voices
Voice-to-voice cloning and musical prosody
Make 3D stuff
Models that generate 3D objects, scenes, radiance fields, textures and multi-views.
Chat with images
Ask language models about images
Extract text from images
Optical character recognition (OCR) and text extraction
Use FLUX fine-tunes
Browse the diverse range of fine-tunes the community has custom-trained on Replicate
Control image generation
Guide image generation with more than just text. Use edge detection, depth maps, and sketches to get the results you want.
Popular models
SDXL-Lightning by ByteDance: a fast text-to-image model that makes high-quality images in 4 steps
This is an optimised version of the FLUX.1 [schnell] model from Black Forest Labs made with Pruna. We achieve a ~3x speedup over the original model with minimal quality loss.
Return CLIP features for the clip-vit-large-patch14 model
Practical face restoration algorithm for *old photos* or *AI-generated faces*
Generate CLIP (clip-vit-large-patch14) text & image embeddings
Latest models
A combination of ip_adapter SDv1.5 and mediapipe-face to inpaint a face
The Yi series models are large language models trained from scratch by developers at 01.AI.
The Yi series models are large language models trained from scratch by developers at 01.AI.
The Yi series models are large language models trained from scratch by developers at 01.AI.
Custom improvements like a custom callback to enhance the inference | It's a WIP and it may causes some wrong outputs
An extremely fast all-in-one model to use LCM with SDXL, ControlNet and custom LoRA url's!
Create variations of an uploaded image. Please see README for more details
Source: meta-llama/Llama-2-7b-chat-hf ✦ Quant: TheBloke/Llama-2-7B-Chat-AWQ ✦ Intended for assistant-like chat
Source: meta-math/MetaMath-Mistral-7B ✦ Quant: TheBloke/MetaMath-Mistral-7B-AWQ ✦ Bootstrap Your Own Mathematical Questions for Large Language Models
Source: Severian/ANIMA-Phi-Neptune-Mistral-7B ✦ Quant: TheBloke/ANIMA-Phi-Neptune-Mistral-7B-AWQ ✦ Biomimicry Enhanced LLM
Animate Your Personalized Text-to-Image Diffusion Models (Long boot times!)
Latent Consistency Model (LCM): SDXL, distills the original model into a version that requires fewer steps (4 to 8 instead of the original 25 to 50)
Generate high-quality images faster with Latent Consistency Models (LCM), a novel approach that distills the original model, reducing the steps required from 25-50 to just 4-8 in Stable Diffusion (SDXL) image generation.
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDXL images with an image prompt
Canny, soft edge, depth, lineart, segmentation, pose, etc
POC of SDXL-LCM LoRA combined with a Replicate LoRA for 4 second inference time
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDXL images with an image prompt
Latent Consistency Model (LCM): SSD-1B, is a LCM distilled version that reduces the number of inference steps needed to only 2 - 8 steps
OpenChat: Advancing Open-source Language Models with Mixed-Quality Data
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDv1.5 images with an image prompt
Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities
Batch mode for Segmind Stable Diffusion Model (SSD-1B) txt2img
ThinkDiffusionXL is a go-to model capable of amazing photorealism that's also versatile enough to generate high-quality images across a variety of styles and subjects without needing to be a prompting genius
